Ah_San
The v1.0 model is mainly a combination of lora and animage3.1-xl.
v1.0模型主要通过lora与animagine3.1-xl的合并。
The main features are the relatively fixed lo1i face shape and the relatively stable nai style.
主要特色是偏固定的lo1i脸型与较稳定的nai系画风。
v1.0 is an attempt, if you want to use it:
v1.0是一次尝试，如果想使用它可以：
Positive words:
masterpiece, best quality, very aesthetic, absurdres, Hint words for starting:
masterpiece, best quality, Negative prompt:
artist name, watermark, signature, lowres, bad anatomy, bad hands, text, error, missing fingers, extra digit, fewer digits, cropped, worst quality, low quality, normal quality, jpeg artifacts, signature, watermark, username, blurry, artist name worst quality, low quality, swollen or deformed toes, overlapping toes, excess toenail, incorrect toe count, uneven toe count, unnatural foot proportions, Hire. fix: R-ESRGAN 4x+ Anime6B, 4x-UltraSharp, 4x-AnimeSharp,High score iteration steps: 10+,Redraw amplitude: 0.3-0.5,What is Ath_LorimagineXL?
Ath_LorimagineXL is a highly specialized Image generation AI Model of type Safetensors / Checkpoint AI Model created by AI community user Ah_San. Derived from the powerful Stable Diffusion (SDXL 1.0) model, Ath_LorimagineXL has undergone an extensive fine-tuning process, leveraging the power of a dataset consisting of images generated by other AI models or user-contributed data. This fine-tuning process ensures that Ath_LorimagineXL is capable of generating images that are highly relevant to the specific use-cases it was designed for, such as anime, base model, girls.
With a rating of 0 and over 0 ratings, Ath_LorimagineXL is a popular choice among users for generating high-quality images from text prompts.
Can I download Ath_LorimagineXL?
Yes! You can download the latest version of Ath_LorimagineXL from here.
How to use Ath_LorimagineXL?
To use Ath_LorimagineXL, download the model checkpoint file and set up an UI for running Stable Diffusion models (for example, AUTOMATIC1111). Then, provide the model with a detailed text prompt to generate an image. Experiment with different prompts and settings to achieve the desired results. If this sounds a bit complicated, check out our initial guide to Stable Diffusion – it might be of help. And if you really want to dive deep into AI image generation and understand how set up AUTOMATIC1111 to use Safetensors / Checkpoint AI Models like Ath_LorimagineXL, check out our crash course in AI image generation.
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Latest version (v1.0): 1 File
About this version: v1.0
v1.0模型主要通过lora与animagine3.1-xl的合并。
主要特色是偏固定的lo1i脸型与较稳定的nai系画风。
The v1.0 model is mainly a combination of lora and animage3.1-xl.
The main features are the relatively fixed lo1i face shape and the relatively stable nai style.
1 Version
Go ahead and upload yours!
