Stable Diffusion is a text-to-image model by StabilityAI.
Stable Diffusion is a powerful artificial intelligence model capable of generating high-quality images based on text descriptions. Developed by Stability AI in collaboration with various academic researchers and non-profit organizations in 2022, it takes a piece of text and creates an image that closely aligns with the description provided. This model can be used in a variety of applications including image creation, image editing, and even image translation based on text prompts.
The underlying technology of Stable Diffusion is a type of deep learning network known as a latent diffusion model. The process starts with compressing the image from pixel space to a smaller dimensional latent space using a component called a Variational Autoencoder (VAE). The model then applies Gaussian noise to the compressed image and uses a U-Net block to clean up this noise and restore the image to its original form. The final image is generated by converting the representation back into pixel space.
What makes Stable Diffusion unique is its ability to be 'conditioned' on a string of text, an image, or another modality. This means that it can generate images based on a given text prompt or alter an existing image according to the prompt. Additionally, unlike its predecessors like DALL-E and Midjourney, Stable Diffusion has made its code and model weights publicly available, which makes it accessible for individual developers and researchers.
Despite its impressive capabilities, Stable Diffusion does have some limitations. It struggles with certain types of images such as human limbs and faces due to insufficient training data, and it requires significant computing power to train on new data. Additionally, it's worth noting that the model was primarily trained on images with English descriptions, which can result in a bias towards western perspectives and cultures.
Despite these challenges, Stable Diffusion represents a significant step forward in the field of text-to-image AI models. It offers a wealth of possibilities for artists, developers, and researchers alike, enabling them to generate and manipulate images in ways that were previously only possible with extensive human effort and expertise.
Stable Diffusion also provides some unique capabilities that are not found in previous text-to-image models like DALL-E and Midjourney. One of these is the use of textual inversions and LoRAs, or "Latent Optimizers over Randomly Initialized Architectures". Textual inversions allow users to create "embeddings" from a collection of their own images, essentially enabling the model to generate images similar to those in the collection whenever specific words or phrases are used in a text prompt. This capability can be used to reduce biases within the original model or to mimic particular visual styles. LoRAs, on the other hand, are a technique that helps guide the model towards specific types of outputs, such as imitating the style of a particular artist.
Another exciting feature of Stable Diffusion is the ability for users to train their own fine-tuned models. With this capability, users can tailor the model to generate images that cater to specific use-cases, creating outputs that are more aligned with their unique needs and preferences. Techniques such as ControlNet and DreamBooth further enhance this capability. ControlNet is a neural network architecture designed to manage diffusion models by incorporating additional conditions, preserving the integrity of the original model while learning new conditions. DreamBooth, on the other hand, is a fine-tuning model that generates precise, personalized outputs depicting a specific subject based on a set of images. These features make Stable Diffusion a highly adaptable tool that can be customized to generate a broad range of image outputs based on text prompts.